Stable Diffusion 1.5

v1.5 refines the SD1 line for cleaner faces, steadier composition, and better prompt adherence while keeping the lightweight latent pipeline that made SD popular. You can generate from text, guide with a reference image, or mask regions for precise edits, then upscale for delivery. The ecosystem is mature: checkpoints, LoRAs, ControlNet conditions, and schedulers make style and layout control predictable, and negative prompts help suppress unwanted artifacts. For teams and creators, v1.5 remains a dependable, self-hosted foundation for concept art, product shots, posters, and UI mockups with repeatable, on-brand results at low cost.

Overview

Stable Diffusion v1.5 is the classic open latent diffusion model for text-to-image. It runs fast on consumer GPUs, supports image-to-image and inpainting, and serves as a strong base for fine-tuning, LoRA adapters, and ControlNet workflows.

🖼️Image generation 🔊Text to speech 🔍Image upscaling 🖌️Image editing

About Stability AI

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Industry: Artificial Intelligence

Company Size: 184

Location: London, GB

Website: stability.ai

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